Animagine XL 2.0
| Property | Value |
|---|---|
| Author | Linaqruf |
| License | CreativeML OpenRAIL++ |
| Base Model | Stable Diffusion XL 1.0 |
| Training Data | 170k + 83k images |
| Architecture | StableDiffusionXLPipeline |
What is Animagine XL 2.0?
Animagine XL 2.0 is an advanced latent text-to-image diffusion model specifically designed for creating high-resolution anime-style images. Built upon Stable Diffusion XL 1.0, it represents a significant upgrade from its predecessor, incorporating extensive training on a dataset of over 250,000 high-quality anime images.
Implementation Details
The model was trained on an A100 80GB GPU using a two-stage approach: a feature alignment stage with 170k images and an aesthetic tuning stage with 83k high-quality synthetic datasets. It implements advanced training configurations including a learning rate of 1e-6, batch size of 32, and mixed-precision fp16 training.
- Supports multiple aspect ratios from 1024x1024 to 1536x640
- Includes specialized LoRA adapters for style customization
- Implements quality tag system for output control
- Features improved pose handling and character consistency
Core Capabilities
- High-quality anime image generation with consistent style
- Advanced pose and perspective handling
- Multiple resolution support for various use cases
- Integrated quality control through tag system
- Specialized LoRA collection for style adaptation
Frequently Asked Questions
Q: What makes this model unique?
The model combines extensive training data, advanced architecture, and specialized anime-focused optimization. Its unique quality tag system and LoRA adapters allow for precise control over output style and quality.
Q: What are the recommended use cases?
The model excels in creating anime-style artwork for entertainment, media production, art and design, and educational purposes. It's particularly effective for character illustrations and scenic compositions in anime style.
